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  2. Stable Diffusion Version 2 - GitHub

    github.com/Stability-AI/StableDiffusion

    Stable Diffusion is an open-source latent diffusion model for high-resolution image synthesis available on GitHub.

  3. AUTOMATIC1111/stable-diffusion-webui - GitHub

    github.com/AUTOMATIC1111/stable-diffusion-webui

    Stable Diffusion web UI. Contribute to AUTOMATIC1111/stable-diffusion-webui development by creating an account on GitHub.

  4. CompVis/stable-diffusion - GitHub

    github.com/CompVis/stable-diffusion

    Stable Diffusion v1 refers to a specific configuration of the model architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet and CLIP ViT-L/14 text encoder for the diffusion model. The model was pretrained on 256x256 images and then finetuned on 512x512 images. Note: Stable Diffusion v1 is a general text-to-image diffusion ...

  5. stable-diffusion/README.md at main - GitHub

    github.com/CompVis/stable-diffusion/blob/main/README.md

    A latent text-to-image diffusion model. Contribute to CompVis/stable-diffusion development by creating an account on GitHub.

  6. Kameronski/stable-diffusion-1.5 - GitHub

    github.com/Kameronski/stable-diffusion-1.5

    Stable Diffusion v1 refers to a specific configuration of the model architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet and CLIP ViT-L/14 text encoder for the diffusion model. The model was pretrained on 256x256 images and then finetuned on 512x512 images. Note: Stable Diffusion v1 is a general text-to-image diffusion ...

  7. Stable Diffusion XL - GitHub

    github.com/huggingface/diffusers/blob/main/docs/source/en/using-diffusers/sdxl.md

    Stable Diffusion XL. [ [open-in-colab]] Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of ...

  8. Stable Diffusion 3 - GitHub

    github.com/.../docs/source/en/api/pipelines/stable_diffusion/stable_diffusion_3.md

    Stable Diffusion 3 (SD3) was proposed in Scaling Rectified Flow Transformers for High-Resolution Image Synthesis by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion ...

  9. GitHub - Diffusion-models/stable-diffusion-webui: Stable...

    github.com/Diffusion-models/stable-diffusion-webui

    Stable Diffusion was made possible thanks to a collaboration with Stability AI and Runway and builds upon our previous work: CVPR '22 Oral which is available on GitHub. PDF at arXiv. Please also visit our Project page. Stable Diffusion is a latent text-to-image diffusion model.

  10. lllyasviel/stable-diffusion-webui-forge - GitHub

    github.com/lllyasviel/stable-diffusion-webui-forge

    Stable Diffusion WebUI Forge. Stable Diffusion WebUI Forge is a platform on top of Stable Diffusion WebUI (based on Gradio ) to make development easier, optimize resource management, speed up inference, and study experimental features. The name "Forge" is inspired from "Minecraft Forge". This project is aimed at becoming SD WebUI's Forge.

  11. GitHub - Panchovix/stable-diffusion-webui-reForge

    github.com/Panchovix/stable-diffusion-webui-reForge

    Stable Diffusion WebUI Forge/reForge is a platform on top of Stable Diffusion WebUI (based on Gradio) to make development easier, optimize resource management, speed up inference, and study experimental features. The name "Forge" is inspired from "Minecraft Forge". This project is aimed at becoming SD WebUI's Forge.